The Stable Diffusion XL (SDXL) model is the official upgrade to the v1. Conclusion This script is a comprehensive example of. Hugging Face. Its superior capabilities, user-friendly interface, and this comprehensive guide make it an invaluable. Why are my SDXL renders coming out looking deep fried? analog photography of a cat in a spacesuit taken inside the cockpit of a stealth fighter jet, fujifilm, kodak portra 400, vintage photography Negative prompt: text, watermark, 3D render, illustration drawing Steps: 20, Sampler: DPM++ 2M SDE Karras, CFG scale: 7, Seed: 2582516941, Size: 1024x1024,. You can disable this in Notebook settings However, SDXL doesn't quite reach the same level of realism. VRAM settings. Just every 1 in 10 renders/prompt I get cartoony picture but w/e. Text-to-Image • Updated 7 days ago • 361 • 2 Nacken/Gen10. 9 Model. Make sure to upgrade diffusers to >= 0. How to use the Prompts for Refine, Base, and General with the new SDXL Model. Optional: Stopping the safety models from. I run on an 8gb card with 16gb of ram and I see 800 seconds PLUS when doing 2k upscales with SDXL, wheras to do the same thing with 1. SDXL is a new checkpoint, but it also introduces a new thing called a refiner. It's beter than a complete reinstall. No warmaps. The SDXL model is a new model currently in training. Include private repos Repository: . 9 produces massively improved image and composition detail over its predecessor. The basic steps are: Select the SDXL 1. This repository provides the simplest tutorial code for developers using ControlNet with. . This is interesting because it only upscales in one step, without having to take it. How Use Stable Diffusion, SDXL, ControlNet, LoRAs For FREE Without A GPU On. py file in it. Yeah SDXL setups are complex as fuuuuk, there are bad custom nodes that do it but the best ways seem to involve some prompt reorganization which is why I do all the funky stuff with the prompt at the start. Overview Load pipelines, models, and schedulers Load and compare different schedulers Load community pipelines and components Load safetensors Load different Stable Diffusion formats Load adapters Push files to the Hub. 1 text-to-image scripts, in the style of SDXL's requirements. While not exactly the same, to simplify understanding, it's basically like upscaling but without making the image any larger. ) Cloud - Kaggle - Free. 0 given by a panel of expert art critics. All prompts share the same seed. Anaconda 的安裝就不多做贅述，記得裝 Python 3. SDXL is great and will only get better with time, but SD 1. Although it is not yet perfect (his own words), you can use it and have fun. 0 to 10. My hardware is Asus ROG Zephyrus G15 GA503RM with 40GB RAM DDR5-4800, two M. ai@gmail. Further development should be done in such a way that Refiner is completely eliminated. 5 model. Loading & Hub. r/StableDiffusion. 0 (SDXL) this past summer. As diffusers doesn't yet support textual inversion for SDXL, we will use cog-sdxl TokenEmbeddingsHandler class. 6 contributors; History: 8 commits. 1 to gather feedback from developers so we can build a robust base to support the extension ecosystem in the long run. to Hilton Head Island). SDXL is a latent diffusion model, where the diffusion operates in a pretrained, learned (and fixed) latent space of an autoencoder. Open the "scripts" folder and make a backup copy of txt2img. Each painting also comes with a numeric score from 0. 0 Depth Vidit, Depth Faid Vidit, Depth, Zeed, Seg, Segmentation, Scribble. 6k hi-res images with randomized prompts, on 39 nodes equipped with RTX 3090 and RTX 4090 GPUs. xls, . What is SDXL model. 0 is released under the CreativeML OpenRAIL++-M License. SDXL 1. ago. main. Although it is not yet perfect (his own words), you can use it and have fun. 🧨 Diffusers Lecture 18: How Use Stable Diffusion, SDXL, ControlNet, LoRAs For FREE Without A GPU On Kaggle Like Google Colab. Aspect Ratio Conditioning. weight: 0 to 5. Using SDXL base model text-to-image. yaml extension, do this for all the ControlNet models you want to use. finally , AUTOMATIC1111 has fixed high VRAM issue in Pre-release version 1. I refuse. Rename the file to match the SD 2. If you want a fully latent upscale, make sure the second sampler after your latent upscale is above 0. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders ( OpenCLIP-ViT/G and CLIP-ViT/L ). In the last few days I've upgraded all my Loras for SD XL to a better configuration with smaller files. hf-import-sdxl-weights Updated 2 months, 4 weeks ago 24 runs sdxl-text Updated 3 months ago 84 runs real-esrgan-a40. Its APIs can change in future. 9 and Stable Diffusion 1. Stable Diffusion XL (SDXL), is the latest AI image generation model that can generate realistic faces, legible text within the images, and better image composition, all while using shorter and simpler prompts. Compare base models. This capability, once restricted to high-end graphics studios, is now accessible to artists, designers, and enthusiasts alike. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. Follow me here by clicking the heart ️ and liking the model 👍, and you will be notified of any future versions I release. safetensors. 49. 既にご存じの方もいらっしゃるかと思いますが、先月Stable Diffusionの最新かつ高性能版である Stable Diffusion XL が発表されて話題になっていました。. 9 espcially if you have an 8gb card. In the Comfyui SDXL workflow example, the refiner is an integral part of the generation process. i git pull and update from extensions every day. 5 for inpainting details. 183. Next as usual and start with param: withwebui --backend diffusers. Further development should be done in such a way that Refiner is completely eliminated. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. He continues to train others will be launched soon. This repository provides the simplest tutorial code for developers using ControlNet with. It's saved as a txt so I could upload it directly to this post. He puts out marvelous Comfyui stuff but with a paid Patreon and Youtube plan. 5 and Steps to 3 Step 4) Generate images in ~<1 second (instantaneously on a 4090) Basic LCM Comfy. The addition of the second model to SDXL 0. xlsx). 🤗 AutoTrain Advanced. It’s designed for professional use, and. Text-to-Image Diffusers ControlNetModel stable-diffusion-xl stable-diffusion-xl-diffusers controlnet. 0-small; controlnet-depth-sdxl-1. Tout d'abord, SDXL 1. This would only be done for safety concerns. DucHaiten-AIart-SDXL; SDXL 1. Stable Diffusion XL has been making waves with its beta with the Stability API the past few months. sdxl_vae. 🧨 DiffusersSD 1. If you've ev. This guide will show you how to use the Stable Diffusion and Stable Diffusion XL (SDXL) pipelines with ONNX Runtime. You'll see that base SDXL 1. Data from Excel spreadsheets (. This score indicates how aesthetically pleasing the painting is - let's call it the 'aesthetic score'. 0 02:52. To run the model, first install the latest version of the Diffusers library as well as peft. Here is the link to Joe Penna's reddit post that you linked to over at Civitai. vae is not necessary with vaefix model. In this quick episode we do a simple workflow where we upload an image into our SDXL graph inside of ComfyUI and add additional noise to produce an altered i. It can produce 380 million gallons of renewable diesel annually. Stable Diffusion AI Art: 1024 x 1024 SDXL image generated using Amazon EC2 Inf2 instance. SDXL has some parameters that SD 1 / 2 didn't for training: original image size: w_original, h_original and crop coordinates: c_top and c_left (where the image was cropped, from the top-left corner) So no more random cropping during training, and no more heads cut off during inference. 335 MB darkside1977 • 2 mo. You want to use Stable Diffusion, use image generative AI models for free, but you can't pay online services or you don't have a strong computer. {"payload":{"allShortcutsEnabled":false,"fileTree":{"torch-neuronx/inference":{"items":[{"name":"customop_mlp","path":"torch-neuronx/inference/customop_mlp. This is just a simple comparison of SDXL1. 5 context, which proves that 1. Description: SDXL is a latent diffusion model for text-to-image synthesis. but when it comes to upscaling and refinement, SD1. main. You can refer to some of the indicators below to achieve the best image quality : Steps : > 50. 9 through Python 3. output device, e. There are 18 high quality and very interesting style Loras that you can use for personal or commercial use. patrickvonplaten HF staff. 0 02:52. Tout d'abord, SDXL 1. Top SDF Flights to International Cities. I don't use --medvram for SD1. He published on HF: SD XL 1. Deepfloyd when it was released few months ago seem to be much better than Midjourney and SD at the time, but need much more Vram. Using the SDXL base model on the txt2img page is no different from using any other models. The first step to using SDXL with AUTOMATIC1111 is to download the SDXL 1. This can usually. 使用 LCM LoRA 4 步完成 SDXL 推理 . LLM: quantisation, fine tuning. "New stable diffusion model (Stable Diffusion 2. An astronaut riding a green horse. And + HF Spaces for you try it for free and unlimited. The SDXL 1. 97 per. There are a few more complex SDXL workflows on this page. 0 model. 0. Stable Diffusion XL (SDXL) 1. 21, 2023. We offer cheap direct, non-stop flights. Today, Stability AI announces SDXL 0. 0. Make sure you go to the page and fill out the research form first, else it won't show up for you to download. Not even talking about. Generated by Finetuned SDXL. T2I Adapter is a network providing additional conditioning to stable diffusion. Although it is not yet perfect (his own words), you can use it and have fun. 5 however takes much longer to get a good initial image. co Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: ; the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters Edit: In case people are misunderstanding my post: This isn't supposed to be a showcase of how good SDXL or DALL-E 3 is at generating the likeness of Harrison Ford or Lara Croft (SD has an endless advantage at that front since you can train your own models), and it isn't supposed to be an argument that one model is overall better than the other. These are the 8 images displayed in a grid: LCM LoRA generations with 1 to 8 steps. He continues to train others will be launched soon!Stable Diffusion XL delivers more photorealistic results and a bit of text. Stability AI. You signed in with another tab or window. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. May need to test if including it improves finer details. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. I haven’t used that particular SDXL openpose model but I needed to update last week to get sdxl controlnet IP-adapter to work properly. Available at HF and Civitai. Tensor values are not checked against, in particular NaN and +/-Inf could be in the file. This repository hosts the TensorRT versions of Stable Diffusion XL 1. Spaces that are too early or cutting edge for mainstream usage 🙂 SDXL ONLY. ComfyUI Impact Pack. Stable Diffusion XL (SDXL) was proposed in SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach. 2-0. 5 on A1111 takes 18 seconds to make a 512x768 image and around 25 more seconds to then hirezfix it to 1. OS= Windows. Also try without negative prompts first. 5: 512x512 SD 1. Stable Diffusion XL (SDXL) was proposed in SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis by Dustin Podell, Zion English, Kyle Lacey, Andreas. Imagine we're teaching an AI model how to create beautiful paintings. Commit. With Automatic1111 and SD Next i only got errors, even with -lowvram parameters, but Comfy. Adetail for face. RENDERING_REPLICATE_API_MODEL: optional, defaults to "stabilityai/sdxl" RENDERING_REPLICATE_API_MODEL_VERSION: optional, in case you want to change the version; Language model config: LLM_HF_INFERENCE_ENDPOINT_URL: "" LLM_HF_INFERENCE_API_MODEL:. And + HF Spaces for you try it for free and unlimited. 0 的过程，包括下载必要的模型以及如何将它们安装到. SDXL, ControlNet, Nodes, in/outpainting, img2img, model merging, upscaling, LORAs,. Install SD. KiwiSDR sound client for Mac by Black Cat Systems. This commit does not belong to any branch on this repository, and may belong to a fork outside of the repository. x with ControlNet, have fun!camenduru/T2I-Adapter-SDXL-hf. The pre-trained models showcase a wide-range of conditions, and the community has built others, such as conditioning on pixelated color palettes. It adds pairs of rank-decomposition weight matrices (called update matrices) to existing weights, and only trains those newly added weights. In the AI world, we can expect it to be better. 9 model , and SDXL-refiner-0. x ControlNet's in Automatic1111, use this attached file. 9 sets a new benchmark by delivering vastly enhanced image quality and. 5GB vram and swapping refiner too , use --medvram-sdxl flag when starting r/StableDiffusion • Year ahead - Requests for Stability AI from community? The disadvantage is that slows down generation of a single image SDXL 1024x1024 by a few seconds for my 3060 GPU. SDXL Inpainting is a desktop application with a useful feature list. Kohya_ss has started to integrate code for SDXL training support in his sdxl branch. Register for your free account. This is my current SDXL 1. This process can be done in hours for as little as a few hundred dollars. gr-kiwisdr GNURadio support for KiwiSDR by. Optionally, we have just added a new theme, Amethyst-Nightfall, (It's purple!) you can select that at the top in UI theme. See the usage instructions for how to run the SDXL pipeline with the ONNX files hosted in this repository. The final test accuracy is 89. Generate comic panels using a LLM + SDXL. Upscale the refiner result or dont use the refiner. Apologies if this has already been posted, but Google is hosting a pretty zippy (and free!) HuggingFace Space for SDXL. Rename the file to match the SD 2. positive: more realistic. In principle you could collect HF from the implicit tree-traversal that happens when you generate N candidate images from a prompt and then pick one to refine. . made by me) requests an image using an SDXL model, they get 2 images back. Replicate SDXL LoRAs are trained with Pivotal Tuning, which combines training a concept via Dreambooth LoRA with training a new token with Textual Inversion. Now go enjoy SD 2. I see that some discussion have happend here #10684, but having a dedicated thread for this would be much better. All the controlnets were up and running. 1 recast. Available at HF and Civitai. Reload to refresh your session. 5 and 2. Although it is not yet perfect (his own words), you can use it and have fun. Most comprehensive LORA training video. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. You can then launch a HuggingFace model, say gpt2, in one line of code: lep photon run --name gpt2 --model hf:gpt2 --local. AutoTrain Advanced: faster and easier training and deployments of state-of-the-art machine learning models. The model learns by looking at thousands of existing paintings. Download the SDXL 1. CFG : 9-10. You can assign the first 20 steps to the base model and delegate the remaining steps to the refiner model. As diffusers doesn't yet support textual inversion for SDXL, we will use cog-sdxl TokenEmbeddingsHandler class. The post just asked for the speed difference between having it on vs off. Aug. - various resolutions to change the aspect ratio (1024x768, 768x1024, also did some testing with 1024x512, 512x1024) - upscaling 2X with Real-ESRGAN. LCM 模型 (Latent Consistency Model) 通过将原始模型蒸馏为另一个需要更少步数 (4 到 8 步，而不是原来的 25 到 50 步) 的版本以减少用 Stable Diffusion (或 SDXL) 生成图像所需的步数。. 0 mixture-of-experts pipeline includes both a base model and a refinement model. We release T2I-Adapter-SDXL, including sketch, canny, and keypoint. ) Stability AI. Updating ControlNet. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. Description for enthusiast AOM3 was created with a focus on improving the nsfw version of AOM2, as mentioned above. InoSim. 9 was meant to add finer details to the generated output of the first stage. You really want to follow a guy named Scott Detweiler. 🤗 Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. 0 enhancements include native 1024-pixel image generation at a variety of aspect ratios. civitAi網站1. Although it is not yet perfect (his own words), you can use it and have fun. Stable Diffusion XL (SDXL) 1. like 387. 22 Jun. It’s important to note that the model is quite large, so ensure you have enough storage space on your device. 0 created in collaboration with NVIDIA. 1 - SDXL UI Support, 8GB VRAM, and More. 1 reply. He published on HF: SD XL 1. 🧨 Diffusers SD 1. 5 model, if using the SD 1. The application isn’t limited to just creating a mask within the application, but extends to generating an image using a text prompt and even storing the history of your previous inpainting work. Describe the image in detail. Stable Diffusion XL (SDXL) is the latest AI image model that can generate realistic people, legible text, and diverse art styles with excellent image composition. I noticed the more bizarre your prompt gets, the more SDXL wants to turn it into a cartoon. SDXL models are really detailed but less creative than 1. 0 VAE, but when I select it in the dropdown menu, it doesn't make any difference (compared to setting the VAE to "None"): images are exactly the same. 0, an open model representing the next evolutionary. ckpt) and trained for 150k steps using a v-objective on the same dataset. SDXL 1. Although it is not yet perfect (his own words), you can use it and have fun. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. 5 and SD v2. The abstract from the paper is: We present SDXL, a latent diffusion model for text-to-image synthesis. Independent U. SDXL is supposedly better at generating text, too, a task that’s historically. 5 LoRA: Link: HF Link: We then need to include the LoRA in our prompt, as we would any other LoRA. DocumentationThe chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. I'm already in the midst of a unique token training experiment. Downscale 8 times to get pixel perfect images (use Nearest Neighbors) Use a fixed VAE to avoid artifacts (0. 下載 WebUI. If you would like to access these models for your research, please apply using one of the following links: SDXL-base-0. r/StableDiffusion. For the base SDXL model you must have both the checkpoint and refiner models. See the official tutorials to learn them one by one. If you want to load a PyTorch model and convert it to the ONNX format on-the-fly, set export=True:Make sure your Controlnet extension is updated in the Extension tab, SDXL support has been expanding the past few updates and there was one just last week. Image To Image SDXL tonyassi Oct 13. A new version of Stability AI’s AI image generator, Stable Diffusion XL (SDXL), has been released. Awesome SDXL LoRAs. jbilcke-hf HF staff commited on Sep 7. SDXL Inpainting is a latent diffusion model developed by the HF Diffusers team. 9 and Stable Diffusion 1. Stable Diffusion: - I run SDXL 1. 9 and Stable Diffusion 1. PixArt-Alpha. Low-Rank Adaptation of Large Language Models (LoRA) is a training method that accelerates the training of large models while consuming less memory. The skilled prompt crafter can break away from the "usual suspects" and draw from the thousands of styles of those artists recognised by SDXL. . clone. How to Do SDXL Training For FREE with Kohya LoRA - Kaggle - NO GPU Required - Pwns Google Colab. I have to believe it's something to trigger words and loras. like 852. 0 given by a panel of expert art critics. LCM LoRA, LCM SDXL, Consistency Decoder LCM LoRA. This workflow uses both models, SDXL1. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. SDNEXT, with diffusors and sequential CPU offloading can run SDXL at 1024x1024 with 1. 5 version) Step 3) Set CFG to ~1. Contribute to dai-ma-tai-nan-le/ai- development by creating an account on. He continues to train. Could not load tags. This checkpoint is a LCM distilled version of stable-diffusion-xl-base-1. made by me). sayak_hf 2 hours ago | prev | next [–] The Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of the Stable Diffusion XL (SDXL), offering a 60% speedup while maintaining high-quality text-to-image generation capabilities. Browse sdxl Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAsSDXL ControlNets 🚀. The model is capable of generating images with complex concepts in various art styles, including photorealism, at quality levels that exceed the best image models available today. As diffusers doesn't yet support textual inversion for SDXL, we will use cog-sdxl TokenEmbeddingsHandler class. 5 to inpaint faces onto a superior image from SDXL often results in a mismatch with the base image. Refer to the documentation to learn more. 0 can achieve many more styles than its predecessors, and "knows" a lot more about each style. 9 espcially if you have an 8gb card. 57967/hf/0925. Additionally, there is a user-friendly GUI option available known as ComfyUI. Astronaut in a jungle, cold color palette, muted colors, detailed, 8k. SDXL works "fine" with just the base model, taking around 2m30s to create a 1024x1024 image (SD1. In addition make sure to install transformers, safetensors, accelerate as well as the invisible watermark: pip install invisible_watermark transformers accelerate safetensors. T2I-Adapter aligns internal knowledge in T2I models with external control signals. 0) stands at the forefront of this evolution. N prompt:[Tutorial] How To Use Stable Diffusion SDXL Locally And Also In Google Colab On Google Colab . Generate text2image "Picture of a futuristic Shiba Inu", with negative prompt "text, watermark" using SDXL base 0. An astronaut riding a green horse. Now you can input prompts in the typing area and press Enter to send prompts to the Discord server. All we know is it is a larger model with more parameters and some undisclosed improvements. June 27th, 2023. The model is released as open-source software. You're asked to pick which image you like better of the two. x ControlNet model with a . In the last few days I've upgraded all my Loras for SD XL to a better configuration with smaller files. 5 and they will tell more or less the same. In fact, it may not even be called the SDXL model when it is released. Tiny-SD, Small-SD, and the SDXL come with strong generation abilities out of the box. It achieves impressive results in both performance and efficiency. 0 (no fine-tuning, no LoRA) 4 times, one for each panel ( prompt source code ) - 25 inference steps. LCM LoRA SDXL. huggingface / blog Public. T2I-Adapter aligns internal knowledge in T2I models with external control signals. yaml extension, do this for all the ControlNet models you want to use. Applications in educational or creative tools. If you fork the project you will be able to modify the code to use the Stable Diffusion technology of your choice (local, open-source, proprietary, your custom HF Space etc). SDXL uses base+refiner, the custom modes use no refiner since it's not specified if it's needed. 2 (1Tb+2Tb), it has a NVidia RTX 3060 with only 6GB of VRAM and a Ryzen 7 6800HS CPU. Select bot-1 to bot-10 channel. x ControlNet model with a . latest Nvidia drivers at time of writing. Two-model workflow is a dead-end development, already now models that train based on SDXL are not compatible with Refiner. Input prompts. Contact us to learn more about fine-tuning stable diffusion for your use. and some features, such as using the refiner step for SDXL or implementing upscaling, haven't been ported over yet. Typically, PyTorch model weights are saved or pickled into a . r/StableDiffusion. Anyways, if you’re using “portrait” in your prompt that’s going to lead to issues if you’re trying to avoid it. 0 is released under the CreativeML OpenRAIL++-M License. He published on HF: SD XL 1. SDXL Support for Inpainting and Outpainting on the Unified Canvas. Additionally, it accurately reproduces hands, which was a flaw in earlier AI-generated images. In the last few days I've upgraded all my Loras for SD XL to a better configuration with smaller files. The following SDXL images were generated on an RTX 4090 at 1280×1024 and upscaled to 1920×1152, in 4. Although it is not yet perfect (his own words), you can use it and have fun. Stable Diffusion XL（通称SDXL）の導入方法と使い方. On some of the SDXL based models on Civitai, they work fine. Now go enjoy SD 2. The SDXL refiner is incompatible and you will have reduced quality output if you try to use the base model refiner with ProtoVision XL. 0 model will be quite different. To just use the base model, you can run: import torch from diffusers import. This ability emerged during the training phase of the AI, and was not programmed by people. We would like to show you a description here but the site won’t allow us. py with model_fn and optionally input_fn, predict_fn, output_fn, or transform_fn. reply. There were any NSFW SDXL models that were on par with some of the best NSFW SD 1. SD. Tasks. Full tutorial for python and git. There are also FAR fewer LORAs for SDXL at the moment. Introduced with SDXL and usually only used with SDXL based models, it's meant to come in at the last x amount of generation steps instead of the main model to add detail to the image. SD-XL. Too scared of a proper comparison eh. . SDXL 0. SargeZT has published the first batch of Controlnet and T2i for XL. Join. UJL123 • 3 mo. 0 with some of the current available custom models on civitai. 9 Research License. safetensor version (it just wont work now) Downloading model. The latent output from step 1 is also fed into img2img using the same prompt, but now using "SDXL_refiner_0.